Tutorials

Local Installation of Stable Diffusion 3.5

So, it’s finally here. Stable Diffusion 3.5 has been released by StabilityAI on October 22nd, 2024. After huge back clashes in the community on Stable Diffusion 3, they are back with the improved version. Basically three model variants are on the…

Video Depth Mapper: Efficient 3D Depth Mapping Solution

Due to the extreme diversity of video content like motion, camera panning, and length, the challenging part arises when working with video frames to attain depth estimation. Depth Crafter can make your life easier. It has been released by Tencent AI Lab,3ARC Lab,…

Top 30 Negative Prompts for Reliable Diffusion Models

  You can use these negative prompts while generating images using Stable Diffusion models. 1. Clear Branding Prompt: text, logos, watermarks 2. Sharp Visuals Prompt: blurry backgrounds, distinct features 3. Anatomical Accuracy Prompt: distorted body parts, anatomical inaccuracies 4. Soft…

19 Captivating Selfie Ideas to Slay Your Feed

The most important problem many people face is how to take a selfie shot. Here are some amazing prompts, we created and tested on Flux, Stable diffusion XL (SDXL), and other Stable Diffusion models.  These will be helpful if you…

ComfyUI: Transform Images/Text into Lengthy Videos with Pyramid Flow

Generating longer video with maximum consistency is one of the challenging task. But now it can be possible with Pyramid Flow. A text to video open source model based on Stable Diffusion3 Medium, CogVideoX, Flux1.0, WebVid-10M, OpenVid-1M, Diffusion Forcing, GameNGen,…

Enhance Image Prompts with TIPO and DanTagGen

Are you searching for the best LLM that can optimize your prompts? Here it is. TIPO and DanTagGen are the LLMs(large language models) that help you to “text-pre-sampling” in the workflow. It generates the detailed tags relevant to your inputted…

Enhance Your Images: Elevate with SUPIR

Forget about those older Photoshop techniques where you need editing skills but don’t get the perfect resolution. SUPIR (Scaling-UP Image Restoration) based on LoRA and Stable Diffusion XL(SDXL) framework released by XPixel group , helps you to upscale your image…

Enhanced ControlNet Alpha for Flux Inpainting in ComfyUI

Inpainting ControlNet Alpha model for the FLUX.1-dev model released by the Alimama Creative Team works under Alibaba. The respective model weights fall under Flux dev non-commercial license.  Its trained on 12 million images from the laion2B dataset and other internal…

Generate AI Videos with ControlNext and SVD V2

The methodology is to use ControlNext model(released by DV labs research) with SVD V2 (by StabilityAI) to create consistent AI videos.  The actual architecture just been cloned the way AnymateAnyone works. The model has been trained on better, higher-quality videos with humans pose…

FLUX: Streamline Installation Process

The new text-to-image diffusion model Flux is destroying all open-source and black box models. This model has been released by Black Forest Labs. Trained with 12 billion parameters based on multimodal and parallel diffusion transformer block architecture. FLUX : Installation is Here !! 😍 Tested…